Stable Diffusion is a generative AI model developed to create unique photorealistic images from text and image prompts. Launched in 2022, it employs a latent diffusion approach, which allows it to operate efficiently on consumer hardware with modest GPU requirements. The model utilizes a variational autoencoder to compress images into a latent space, enabling it to generate and manipulate high-quality visuals. It supports various functionalities, including text-to-image generation, image inpainting, and outpainting, making it versatile for creative tasks. Stable Diffusion's open-source nature allows users to fine-tune the model with minimal data, enhancing its adaptability for specific needs. As of 2024, Stable Diffusion 3 has introduced advanced prompt-following capabilities, including better understanding of complex prompts with multiple subjects and detailed text rendering. The model now employs a new architecture based on a stack of Diffusion Transformers. Additionally, updated sampling mechanisms, like Rectified Flow sampling, enable faster, higher-quality image generation. Stable Diffusion 3 is available in multiple model sizes, ranging from 800 million to 8 billion parameters, catering to different user needs and hardware capabilities.